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  2. Stable Diffusion - Wikipedia

    en.wikipedia.org/wiki/Stable_Diffusion

    The model is insufficiently trained to understand human limbs and faces due to the lack of representative features in the database, and prompting the model to generate images of such type can confound the model. [37] Stable Diffusion XL (SDXL) version 1.0, released in July 2023, introduced native 1024x1024 resolution and improved generation for ...

  3. ComfyUI - Wikipedia

    en.wikipedia.org/wiki/ComfyUI

    ComfyUI is an open source, node-based program that allows users to generate images from a series of text prompts.It uses free diffusion models such as Stable Diffusion as the base model for its image capabilities combined with other tools such as ControlNet and LCM Low-rank adaptation with each tool being represented by a node in the program.

  4. Fooocus - Wikipedia

    en.wikipedia.org/wiki/Fooocus

    Fooocus is an open source generative artificial intelligence program that allows users to generate images from a text prompt. [3] [4] It uses Stable Diffusion as the base model for its image capabilities as well as a collection of default settings and prompts to make the image generation process more streamlined.

  5. Automatic1111 - Wikipedia

    en.wikipedia.org/wiki/Automatic1111

    AUTOMATIC1111 Stable Diffusion Web UI (SD WebUI, A1111, or Automatic1111 [3]) is an open source generative artificial intelligence program that allows users to generate images from a text prompt. [4] It uses Stable Diffusion as the base model for its image capabilities together with a large set of extensions and features to customize its output.

  6. DreamBooth - Wikipedia

    en.wikipedia.org/wiki/DreamBooth

    The methodology used to run implementations of DreamBooth involves the fine-tuning the full UNet component of the diffusion model using a few images (usually 3--5) depicting a specific subject. Images are paired with text prompts that contain the name of the class the subject belongs to, plus a unique identifier.

  7. Flux (text-to-image model) - Wikipedia

    en.wikipedia.org/wiki/Flux_(text-to-image_model)

    An improved flagship model, Flux 1.1 Pro was released on 2 October 2024. [ 27 ] [ 28 ] Two additional modes were added on 6 November, Ultra which can generate image at four times higher resolution and up to 4 megapixel without affecting generation speed and Raw which can generate hyper-realistic image in the style of candid photography .

  8. 22 health care predictions for 2025 from medical researchers

    www.aol.com/22-health-care-predictions-2025...

    Researchers from Mass General Brigham, a health care system in Boston, Massachusetts, shared with Fox News Digital some of the scientific developments and breakthroughs they expect to see in 2025.

  9. Diffusion model - Wikipedia

    en.wikipedia.org/wiki/Diffusion_model

    GLIDE (2022-03) [62] is a 3.5-billion diffusion model, and a small version was released publicly. [6] Soon after, DALL-E 2 was released (2022-04). [63] DALL-E 2 is a 3.5-billion cascaded diffusion model that generates images from text by "inverting the CLIP image encoder", the technique which they termed "unCLIP".